Stable diffusion
Stable Diffusion image generation models. Use for image AI.From its SKILL.md
npx -y skills add G1Joshi/Agent-Skills --skill stable-diffusionAssembled from the repository path, not quoted from the project. Check it against their README if it does not work.
One thing to look at
- 12 stars12 stars. Stars are a popularity signal and not a quality one, but at this level it is likely that nobody has read this closely except its author, and you would be relying on your own review.
SKILL.md
1.2 KB, 279 tokens by cl100k_base, as published. Nobody here has run it
Stable Diffusion
Stable Diffusion (by Stability AI) is the open standard for image generation. SD 3.5 (2025) improves prompt adherence and typography.
When to Use
- Control: You need exact composition control (ControlNet).
- Local Generation: Run on your own GPU. No censorship/cost.
- Fine-Tuning: training LoRAs on your own face/product.
Core Concepts
Diffusers Library
The Hugging Face library to run SD pipelines in Python.
LoRA (Low-Rank Adaptation)
Small adapter files (100MB) that add a style or character to the base model.
ComfyUI
The node-based GUI for building complex SD workflows (see separate skill).
Best Practices (2025)
Do:
- Use SD 3.5 Large: For best text rendering.
- Use Flux: The community has largely moved to Flux.1 (Black Forest Labs) alongside SD.
- Use ControlNet: To force the image to follow a specific pose or edge map.
Don't:
- Don't use SD 1.5: Unless you need specific legacy LoRAs. SDXL/SD3/Flux are superior.
References
What ships with it
Read from the repository
Just SKILL.md. No reference files, no scripts.